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Stable Diffusion Prompt Lab

Professional Stable Diffusion prompt generator. Manage positive and negative prompts with weight syntax.

highly detailed, sharp focus
blurry, low quality, low resolution, pixelated, jpeg artifacts, noise, grainy, out of focus, poorly drawn, bad anatomy
Generation Settings
Steps: 30, Sampler: DPM++ 2M Karras, CFG scale: 7
Tips:
  • • Use (word) for emphasis, (word:1.5) for weight
  • • CFG 7-11 is usually optimal for most images
  • • More steps = more detail, but diminishing returns after 40-50

How Stable Diffusion Prompt Lab Works

A Stable Diffusion Prompt Generator is an AI engineering utility used to craft precise instructions for generative models. This tool is essential for AI artists, prompt engineers, and game developers optimizing prompts for Midjourney/SDXL, managing negative prompt lists, and increasing image fidelity.

The "Prompt Engineering" Suite

AI models are picky syntax. This tool helps you:

  1. Style Selector: Apply presets like "Cinematic," "Anime," or "Photorealistic."
  2. Syntax Builder: Automatically formats weights (text:1.1) and Lora triggers <lora:details:1>.
  3. Negative Library: One-click insertion of standard quality boosters (e.g., "worst quality, lowres, watermark").

Frequently Asked Questions

Stable Diffusion is an AI model that turns text into images. A Prompt is the text description you give it (e.g., "cyberpunk cat"). The better the prompt, the better the image.

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